stable-diffusion
Diffusion model
A text-to-image model trained on images and text prompts using a diffusion process
A latent text-to-image diffusion model
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Language: Jupyter Notebook
last commit: 8 months ago
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stevinz/awesome-game-engine-devRelated projects:
| Repository | Description | Stars |
|---|---|---|
 stability-ai/stablediffusion | A software project that enables high-resolution image synthesis through a specific type of generative model using latent diffusion processes. | 39,501 |
| stability-ai/stablecascade | An image generation model that balances efficiency and quality, utilizing a cascade architecture to compress images before training a text-conditional model in a highly compressed latent space. | 6,560 |
 automatic1111/stable-diffusion-webui | A web interface for Stable Diffusion, providing an intuitive and interactive experience for generating images through various models and parameters. | 144,342 |
 ashawkey/stable-dreamfusion | Generates 3D content from text using a combination of neural networks and image synthesis. | 8,351 |
 ai-forever/kandinsky-2 | A multilingual text2image latent diffusion model with improved aesthetics and controllability | 2,774 |
 nateraw/stable-diffusion-videos | Generates videos by manipulating the latent space of Stable Diffusion to create morphing sequences from text prompts. | 4,470 |
 basujindal/stable-diffusion | An optimized version of Stable Diffusion with reduced VRAM usage, sacrificing inference speed. | 3,134 |
 lllyasviel/controlnet | An implementation of a neural network structure to control diffusion models by adding extra conditions. | 30,944 |
 modelscope/diffsynth-studio | A software framework for training and utilizing various types of diffusion models. | 6,641 |
 lykosai/stabilitymatrix | A multi-platform package manager and inference UI for Stable Diffusion, supporting various AI models and tools | 4,995 |
 openai/guided-diffusion | This project is a software implementation of a diffusion model architecture, allowing users to generate synthetic images based on a learned distribution. | 6,366 |
 deforum/stable-diffusion | A latent text-to-image diffusion model that generates high-resolution images from text prompts. | 796 |
 huggingface/diffusers | A PyTorch-based library for training and using state-of-the-art diffusion models to generate images, audio, and 3D structures | 26,676 |
 lucidrains/dalle2-pytorch | An implementation of DALL-E 2's text-to-image synthesis neural network in PyTorch | 11,184 |
 deep-floyd/if | A text-to-image synthesis model with a modular design, utilizing a frozen text encoder and cascaded pixel diffusion modules to generate photorealistic images. | 7,699 |